Trong đó các thành phần và các dữ liệu đã được code lại sao cho tối ưu nhất và đem lại trải nghiệm sử. 全体の流れは以下の通りです。. 3D-controlled video generation with live previews. Classifier guidance is a recently introduced method to trade off mode coverage and sample fidelity in conditional diffusion models post training, in the same spirit as low temperature sampling or truncation in other types of generative models. Dreambooth is considered more powerful because it fine-tunes the weight of the whole model. k. photo of perfect green apple with stem, water droplets, dramatic lighting. waifu-diffusion-v1-4 / vae / kl-f8-anime2. However, I still recommend that you disable the built-in. 0 license Activity. ckpt to use the v1. Microsoft's machine learning optimization toolchain doubled Arc. card classic compact. 学術的な研究結果はほぼ含まない、まさに無知なる利用者の肌感なので、その程度のご理解で. LAION-5B is the largest, freely accessible multi-modal dataset that currently exists. Hot New Top Rising. 663 upvotes · 25 comments. AI. 0 significantly improves the realism of faces and also greatly increases the good image rate. This checkpoint is a conversion of the original checkpoint into. 0. You signed in with another tab or window. Generate the image. ) 不同的采样器在不同的step下产生的效果. Then under the setting Quicksettings list add sd_vae after sd_model_checkpoint. This is no longer the case. Adding Conditional Control to Text-to-Image Diffusion Models (ControlNet) by Lvmin Zhang and Maneesh Agrawala. 2 days ago · Stable Diffusion For Aerial Object Detection. Introduction. This example is based on the training example in the original ControlNet repository. Browse gay Stable Diffusion models, checkpoints, hypernetworks, textual inversions, embeddings, Aesthetic Gradients, and LORAsisketch93 commented Feb 16, 2023. 5 and 2. Step 1: Download the latest version of Python from the official website. Readme License. 今回の動画ではStable Diffusion web UIを用いて、美魔女と呼ばれるようなおばさん（熟女）やおじさんを生成する方法について解説していきます. 1 Release. In this survey, we provide an overview of the rapidly expanding body of work on diffusion models, categorizing the research into three key. 9, the full version of SDXL has been improved to be the world's best open image generation model. Edit model card Update. Step 1: Download the latest version of Python from the official website. Then, download and set up the webUI from Automatic1111. 0 release includes robust text-to-image models trained using a brand new text encoder (OpenCLIP), developed by LAION with support from Stability AI, which greatly improves the quality of the generated images compared to earlier V1 releases. kind of cute? 😅 A bit of detail with a cartoony feel, it keeps getting better! With your support, Too. An optimized development notebook using the HuggingFace diffusers library. Height. Start Creating. You can go lower than 0. You can use it to edit existing images or create new ones from scratch. HCP-Diffusion is a toolbox for Stable Diffusion models based on 🤗 Diffusers. AutoV2. Image. Click on Command Prompt. License: refers to the. They both start with a base model like Stable Diffusion v1. 🖊️ marks content that requires sign-up or account creation for a third party service outside GitHub. そのままでも使用に問題はありませんが、Civitaiのデータをより使いやすくしてくれる拡張機能が「Civitai Helper」です。. It is more user-friendly. You signed out in another tab or window. It's free to use, no registration required. Background. • 5 mo. Stable Diffusion creates an image by starting with a canvas full of noise and denoise it gradually to reach the final output. It bundles Stable Diffusion along with commonly-used features (like SDXL, ControlNet, LoRA, Embeddings, GFPGAN, RealESRGAN, k-samplers, custom VAE etc). Below are some of the key features: – User-friendly interface, easy to use right in the browser. face-swap stable-diffusion sd-webui roop Resources. Originally Posted to Hugging Face and shared here with permission from Stability AI. Wed, November 22, 2023, 5:55 AM EST · 2 min read. Using VAEs. We tested 45 different GPUs in total — everything that has. Model checkpoints were publicly released at the end of August 2022 by. ·. Stars. 1 day ago · Product. Explore Countless Inspirations for AI Images and Art. この記事を読んでいただければ、好きなモデルがきっとみつかるはずです。. Stable Diffusion creator Stability AI has announced that users can now test a new generative AI that animates a single image generated from a text prompt to create. yml file to stable-diffusion-webuiextensionssdweb-easy-prompt-selector ags, and you can add, change, and delete freely. Stable Diffusion v1-5 NSFW REALISM Model Card Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. To run tests using a specific torch device, set RIFFUSION_TEST_DEVICE. Supported use cases: Advertising and marketing, media and entertainment, gaming and metaverse. StableStudio marks a fresh chapter for our imaging pipeline and showcases Stability AI's dedication to advancing open-source development within the AI ecosystem. 10 and Git installed. The above tool is a Stable Diffusion Image Variations model that has been fine-tuned to take multiple CLIP image embeddings as inputs, allowing users to combine the image embeddings from multiple images to mix their concepts and add text concepts for greater variation. Diffusion models have emerged as a powerful new family of deep generative models with record-breaking performance in many applications, including image synthesis, video generation, and molecule design. You should use this between 0. this is the original text tranlsated ->. OK perhaps I need to give an upscale example so that it can be really called "tile" and prove that it is not off topic. Here's a list of the most popular Stable Diffusion checkpoint models . I provide you with an updated tool of v1. 3D-controlled video generation with live previews. 5 or XL. Expand the Batch Face Swap tab in the lower left corner. Whilst the then popular Waifu Diffusion was trained on SD + 300k anime images, NAI was trained on millions. In this post, you will see images with diverse styles generated with Stable Diffusion 1. 5 and 1 weight, depending on your preference. Please use the VAE that I uploaded in this repository. 0. 🎨 Limitless Possibilities: From breathtaking landscapes to futuristic cityscapes, our AI can conjure an array of visuals that match your wildest concepts. I) Main use cases of stable diffusion There are a lot of options of how to use stable diffusion, but here are the four main use cases:. Counterfeit-V3 (which has 2. ckpt -> Anything-V3. ai and search for NSFW ones depending on the style I want (anime, realism) and go from there. 1. 在 models/Lora 目录下，存放一张与 Lora 同名的 . py is ran with. 1-base, HuggingFace) at 512x512 resolution, both based on the same number of parameters and architecture as 2. As with all things Stable Diffusion, the checkpoint model you use will have the biggest impact on your results. Compared with previous numerical PF-ODE solvers such as DDIM, DPM-Solver, LCM-LoRA can be viewed as a plug-in neural PF-ODE solver. Stable Video Diffusion está disponible en una versión limitada para investigadores. It's default ability generated image from text, but the mo. Join. 主にautomatic1111で使う用になっていますが、括弧を書き換えればNovelAI記法にもなると思います。. In the examples I Use hires. 转载自互联网, 视频播放量 328、弹幕量 0、点赞数 6、投硬币枚数 0、收藏人数 1、转发人数 0, 视频作者 上边的真精彩, 作者简介 音乐反应点评，相关视频：【mamamoo】她拒绝所有人，【mamamoo】我甚至没有生气，只是有点恼火。. Use your browser to go to the Stable Diffusion Online site and click the button that says Get started for free. Runtime errorHeavenOrangeMix. The goal of this article is to get you up to speed on stable diffusion. 295 upvotes ·. It’s worth noting that in order to run Stable Diffusion on your PC, you need to have a compatible GPU installed. Fooocus. 如果想要修改. ; Install the lastest version of stable-diffusion-webui and install SadTalker via extension. Civitaiに投稿されているLoraのリンク集です。 アニメ系の衣装やシチュエーションのLoraを中心にまとめてます。 注意事項 雑多まとめなので、効果的なモデルがバラバラな可能性があります キャラクター系Lora、リアル系Lora、画風系Loraは含みません（リアル系は2D絵の報告があれば載せます. ジャンル→内容→prompt. like 9. Enter a prompt, and click generate. 662 forks Report repository Releases 2. It has evolved from sd-webui-faceswap and some part of sd-webui-roop. This comes with a significant loss in the range. Step 1: Go to DiffusionBee’s download page and download the installer for MacOS – Apple Silicon. 295,277 Members. "Diffusion" works by training an artificial neural network to reverse a process of adding "noise" (random pixels) to an image. The notebooks contain end-to-end examples of usage of prompt-to-prompt on top of Latent Diffusion and Stable Diffusion respectively. 画像生成AI (Stable Diffusion Web UI、にじジャーニーなど)で画質を調整するする方法を紹介します。. Thank you so much for watching and don't forg. Style. Stable Diffusion XL (SDXL), is the latest AI image generation model that can generate realistic faces, legible text within the images, and better image composition, all while using shorter and simpler prompts. AI動画用のフォルダを作成する. Learn more. 5 as w. download history blame contribute delete. You will learn the main use cases, how stable diffusion works, debugging options, how to use it to your advantage and how to extend it. 5 Resources →. py --prompt "a photograph of an astronaut riding a horse" --plms. 开启后，只需要点击对应的按钮，会自动将提示词输入到文生图的内容栏。. Although some of that boost was thanks to good old. But what is big news is when a major name like Stable Diffusion enters. We're going to create a folder named "stable-diffusion" using the command line. ControlNet and OpenPose form a harmonious duo within Stable Diffusion, simplifying character animation. 1:7860" or "localhost:7860" into the address bar, and hit Enter. この記事で. In this article we'll feature anime artists that you can use in Stable Diffusion models (NAI Diffusion, Anything V3) as well as the official NovelAI and Midjourney's Niji Mode to get better results. UPDATE DETAIL (中文更新说明在下面) Hello everyone, this is Ghost_Shell, the creator. Stable Diffusion WebUI. It is recommended to use the checkpoint with Stable Diffusion v1-5 as the checkpoint has been trained on it. © Civitai 2023. stable-diffusion. 首先保证自己有一台搭载了gtx 1060及其以上品级显卡的电脑（只支持n卡）下载程序主体，B站很多up做了整合包，这里推荐一个，非常感谢up主独立研究员-星空BV1dT411T7Tz这样就能用SD原始的模型作画了然后在这里下载yiffy. Append a word or phrase with -or +, or a weight between 0 and 2 (1=default), to decrease. The name Aurora, which means 'dawn' in Latin, represents the idea of a new beginning and a fresh start. . My AI received one of the lowest scores among the 10 systems covered in Common Sense’s report, which warns that the chatbot is willing to chat with teen users about sex and alcohol and that it. pickle. Example: set COMMANDLINE_ARGS=--ckpt a. This specific type of diffusion model was proposed in. This file is stored with Git LFS . The Stability AI team is proud to release as an open model SDXL 1. Per default, the attention operation. 6 here or on the Microsoft Store. Microsoft's machine learning optimization toolchain doubled Arc. The t-shirt and face were created separately with the method and recombined. You will see the exact keyword applied to two classes of images: (1) a portrait and (2) a scene. At the time of release (October 2022), it was a massive improvement over other anime models. With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. Hot New Top. Stability AI is thrilled to announce StableStudio, the open-source release of our premiere text-to-image consumer application DreamStudio. Stable Diffusion XL (SDXL) is a powerful text-to-image generation model that iterates on the previous Stable Diffusion models in three key ways: the UNet is 3x larger and SDXL combines a second text encoder (OpenCLIP ViT-bigG/14) with the original text encoder to significantly increase the number of parameters. 1. Stable Diffusion demo. CivitAI is great but it has some issue recently, I was wondering if there was another place online to download (or upload) LoRa files. Install Python on your PC. Edited in AfterEffects. Organize machine learning experiments and monitor training progress from mobile. SDK for interacting with stability. Model checkpoints were publicly released at the end of August 2022 by a collaboration of Stability AI, CompVis, and Runway with support from EleutherAI and LAION. Here are a few things that I generally do to avoid such imagery: I avoid using the term "girl" or "boy" in the positive prompt and instead opt for "woman" or "man". 📘English document 📘中文文档. 万叶真的是太帅了！ 视频播放量 309、弹幕量 0、点赞数 3、投硬币枚数 0、收藏人数 0、转发人数 2, 视频作者 鹤秋幽夜, 作者简介 太阳之所以耀眼，是因为它连尘埃都能照亮，相关视频：枫原万叶，芙宁娜与风伤万叶不同配队测试，枫丹最强阵容白万芙特！白万芙特输出手法！Sensitive Content. Think about how a viral tweet or Facebook post spreads—it's not random, but follows certain patterns. Credit Cost. Detailed feature showcase with images: Original txt2img and img2img modes; One click install and run script (but you still must install python and git) Outpainting; Inpainting; Color Sketch; Prompt Matrix; Stable Diffusion UpscaleMidjourney (v4) Stable Diffusion (DreamShaper) Portraits Content Filter. 24 watching Forks. ; Prompt: SD v1. . 1. png 文件然后 refresh 即可。. Other upscalers like Lanczos or Anime6B tends to smoothen them out, removing the pastel-like brushwork. You can create your own model with a unique style if you want. Use the following size settings to. Download the LoRA contrast fix. Upload vae-ft-mse-840000-ema-pruned. Stable Diffusion is a popular generative AI tool for creating realistic images for various uses cases. You can process either 1 image at a time by uploading your image at the top of the page. This is a merge of Pixar Style Model with my own Loras to create a generic 3d looking western cartoon. With Stable Diffusion, you can create stunning AI-generated images on a consumer-grade PC with a GPU. It is a text-to-image generative AI model designed to produce images matching input text prompts. 17 May. 1: SDXL ; 1: Stunning sunset over a futuristic city, with towering skyscrapers and flying vehicles, golden hour lighting and dramatic clouds, high detail, moody atmosphereAnnotated PyTorch Paper Implementations. 第一次做这个，不敢说是教程，分享一下制作的过程，希望能帮到有需要的人, 视频播放量 25954、弹幕量 0、点赞数 197、投硬币枚数 61、收藏人数 861、转发人数 78, 视频作者 ruic-v, 作者简介 ，相关视频：快速把自拍照片动漫化，完全免费！，又来了 ！她带着东西又来了，stable diffusion图生图（真人转. Search. If you click the Option s icon in the prompt box, you can go a little deeper: For Style, you can choose between Anime, Photographic, Digital Art, Comic Book. 免费在线NovelAi智能绘画网站，手机也能用的NovelAI绘画（免费），【Stable Diffusion】在线使用SD 无需部署 无需显卡，在手机上使用stable diffusion，完全免费！. stable-diffusion. This model is a simple merge of 60% Corneo's 7th Heaven Mix and 40% Abyss Orange Mix 3. Stable Diffusion v1. LoRA is added to the prompt by putting the following text into any location: <lora:filename:multiplier> , where filename is the name of file with LoRA on disk, excluding extension, and multiplier is a number, generally from 0 to 1, that lets you choose how. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. In contrast to FP32, and as the number 16 suggests, a number represented by FP16 format is called a half-precision floating point number. 5. . If you can find a better setting for this model, then good for you lol. Stable Diffusion is an AI model launched publicly by Stability. 5, 99% of all NSFW models are made for this specific stable diffusion version. like 9. Part 1: Getting Started: Overview and Installation. {"message":"API rate limit exceeded for 52. Stable Diffusion （ステイブル・ディフュージョン）は、2022年に公開された ディープラーニング （深層学習）の text-to-imageモデル （ 英語版 ） である。. Within this folder, perform a comprehensive deletion of the entire directory associated with Stable Diffusion. 【Termux+QEMU】，手机云端安装运行stable-diffusion-webui教程，【Stable Diffusion】搭建远程AI绘画服务-随时随地用自己的显卡画图，让ChatGPT玩生成艺术？来看看得到了什么~，最大方的AI绘图软件，每天免费画1000张图！【Playground AI绘画教学】. like 66. 7X in AI image generator Stable Diffusion. Low level shot, eye level shot, high angle shot, hip level shot, knee, ground, overhead, shoulder, etc. それでは実際の操作方法について解説します。. The text-to-image fine-tuning script is experimental. You can join our dedicated community for Stable Diffusion here, where we have areas for developers, creatives, and just anyone inspired by this. How to install Stable Diffusion locally ? First, get the SDXL base model and refiner from Stability AI. Stable Diffusion system requirements – Hardware. 小白失踪几天了！. You'll see this on the txt2img tab: An advantage of using Stable Diffusion is that you have total control of the model. Install the Dynamic Thresholding extension. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. Stable Diffusion v2 are two official Stable Diffusion models. In my tests at 512,768 resolution, the good image rate of the Prompts I used before was above 50%. 2023年5月15日 02:52. Stars. The output is a 640x640 image and it can be run locally or on Lambda GPU. Rename the model like so: Anything-V3. Type cmd. In the context of stable diffusion and the current implementation of Dreambooth, regularization images are used to encourage the model to make smooth, predictable predictions, and to improve the quality and consistency of the output images, respectively. Extend beyond just text-to-image prompting. The text-to-image models are trained with a new text encoder (OpenCLIP) and they're able to output 512x512 and 768x768 images. Playing with Stable Diffusion and inspecting the internal architecture of the models. Learn more about GitHub Sponsors. Hi! I just installed the extension following the steps on the readme page, downloaded the pre-extracted models (but the same issue appeared with full models upon trying) and excitedly tried to generate a couple of images, only to see the. I used two different yet similar prompts and did 4 A/B studies with each prompt. 10. Creating Fantasy Shields from a Sketch: Powered by Photoshop and Stable Diffusion. You will find easy-to-follow tutorials and workflows on this site to teach you everything you need to know about. If you don't have the VAE toggle: in the WebUI click on Settings tab > User Interface subtab. Reload to refresh your session. For a minimum, we recommend looking at 8-10 GB Nvidia models. Stable Diffusion XL is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, cultivates autonomous freedom to produce incredible imagery, empowers billions of people to create stunning art within seconds. 顶级AI绘画神器！. Browse logo Stable Diffusion models, checkpoints, hypernetworks, textual inversions, embeddings, Aesthetic Gradients, and LORAs26 Jul. *PICK* (Updated Sep. You switched. An extension of stable-diffusion-webui. Start with installation & basics, then explore advanced techniques to become an expert. 你需要准备同样角度的其他背景色底图用于ControlNet勾线3. The solution offers an industry leading WebUI, supports terminal use through a CLI, and serves as the foundation for multiple commercial products. This is how others see you. 0 的过程，包括下载必要的模型以及如何将它们安装到. 5. Perfect for artists, designers, and anyone who wants to create stunning visuals without any. Linter: ruff Formatter: black Type checker: mypy These are configured in pyproject. The extension is fully compatible with webui version 1. 4c4f051 about 1 year ago. Put WildCards in to extensionssd-dynamic-promptswildcards folder. このコラムは筆者がstablediffsionを使っていくうちに感じた肌感を同じ利用者について「ちょっとこんなんだと思うんだけど？. If you click the Option s icon in the prompt box, you can go a little deeper: For Style, you can choose between Anime, Photographic, Digital Art, Comic Book. Naturally, a question that keeps cropping up is how to install Stable Diffusion on Windows. trained with chilloutmix checkpoints. The first step to getting Stable Diffusion up and running is to install Python on your PC. LMS is one of the fastest at generating images and only needs a 20-25 step count. Time. 很简单！ 方法一. The Stable Diffusion 2. stage 2：キーフレームの画像を抽出. ckpt uses the model a. You'll see this on the txt2img tab: If you've used Stable Diffusion before, these settings will be familiar to you, but here is a brief overview of what the most important options mean:Intel's latest Arc Alchemist drivers feature a performance boost of 2. License. Create better prompts. Live Chat. Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from CompVis, Stability AI and LAION. 鳳えむ (プロジェクトセカイ) straight-cut bangs, light pink hair,bob cut, shining pink eyes, The girl who puts on a pink cardigan on the gray sailor uniform,white collar, gray skirt, In front of cardigan is open, Ootori-Emu, Cheerful smile, フリスク (undertale) undertale,Frisk. We’re on a journey to advance and democratize artificial intelligence through open source and open science. Note: If you want to process an image to create the auxiliary conditioning, external dependencies are required as shown below:Steps. I literally had to manually crop each images in this one and it sucks. Stable Video Diffusion is released in the form of two image-to-video models, capable of generating 14 and 25 frames at customizable frame rates between 3. Our model uses shorter prompts and generates descriptive images with enhanced composition and realistic aesthetics. . It is too big to display, but you can still download it. g. 34k. 在Stable Diffusion软件中，使用ControlNet+模型实现固定物体批量替换背景出图的流程。一、准备好图片：1. You can join our dedicated community for Stable Diffusion here, where we have areas for developers, creatives, and just anyone inspired by this. ControlNet. 0 release includes robust text-to-image models trained using a brand new text encoder (OpenCLIP), developed by LAION with support. 1️⃣ Input your usual Prompts & Settings. (with < 300 lines of codes!) (Open in Colab) Build a Diffusion model (with UNet + cross attention) and train it to generate MNIST images based on the "text prompt". RePaint: Inpainting using Denoising Diffusion Probabilistic Models. You signed out in another tab or window. com Stable Diffusion is a deep learning, text-to-image model released in 2022 based on diffusion techniques. PromptArt. In Stable Diffusion, although negative prompts may not be as crucial as prompts, they can help prevent the generation of strange images. 一口气学完【12种】Multi-controlnet高阶组合用法合集&SD其他最新插件【持续更新】，Stable Diffusion 控制网络ControlNet的介绍和基础使用 全流程教程（教程合集、持续更新），卷破天际！Stable Diffusion-Controlnet-color线稿精准上色之线稿变为商用成品Training process. Text-to-Image with Stable Diffusion. It is our fastest API, matching the speed of its predecessor, while providing higher quality image generations at 512x512 resolution. License: creativeml-openrail-m. Stable Diffusion is a generative artificial intelligence (generative AI) model that produces unique photorealistic images from text and image prompts. 7B6DAC07D7. Also using body parts and "level shot" helps. 6 API acts as a replacement for Stable Diffusion 1. 2. novelai用了下，故意挑了些涩图tag，效果还可以 基于stable diffusion，操作和sd类似 他们的介绍文档 价格主要是订阅那一下有点贵，要10刀，送1000token 一张图5token(512*768)，细化什么的额外消耗token 这方面倒还好，就是买算力了… 充值token 10刀10000左右，其实还行Use Stable Diffusion outpainting to easily complete images and photos online. This step downloads the Stable Diffusion software (AUTOMATIC1111). 画像生成AIであるStable Diffusionは Mage や DreamStudio などを通して、Webブラウザで簡単に利用することも可能です。. 8 hours ago · The makers of the Stable Diffusion tool "ComfyUI" have added support for Stable AI's Stable Video Diffusion models in a new update. Stable Diffusion WebUI Stable Diffusion WebUI is a browser interface for Stable Diffusion, an AI model that can generate images from text prompts or modify existing images with text prompts. 6. 24 watching Forks. Developed by: Stability AI. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. Characters rendered with the model: Cars and Animals. 注：checkpoints 同理~ 方法二. It can be used in combination with Stable Diffusion, such as runwayml/stable-diffusion-v1-5. Our powerful AI image completer allows you to expand your pictures beyond their original borders. This toolbox supports Colossal-AI, which can significantly reduce GPU memory usage. Sep 15, 2022, 5:30 AM PDT. According to the Stable Diffusion team, it cost them around $600,000 to train a Stable Diffusion v2 base model in 150,000 hours on 256 A100 GPUs. It originally launched in 2022. Stable Diffusion v1 refers to a specific configuration of the model architecture that uses a downsampling-factor 8 autoencoder with an 860M UNet and CLIP ViT-L/14 text encoder for the diffusion model. 0 and fine-tuned on 2. card. to make matters even more confusing, there is a number called a token in the upper right. The text-to-image models in this release can generate images with default. joho. It is trained on 512x512 images from a subset of the LAION-5B database. 1-v, HuggingFace) at 768x768 resolution and (Stable Diffusion 2. シート見るのも嫌な人はマスタ化してるものを適当に整形したのを下に貼っておきます。. Additional training is achieved by training a base model with an additional dataset you are. Step 3: Clone web-ui. Although no detailed information is available on the exact origin of Stable Diffusion, it is known that it was trained with millions of captioned images. 5 base model. Stable Diffusion은 독일 뮌헨 대학교 Machine Vision & Learning Group (CompVis) 연구실의 "잠재 확산 모델을 이용한 고해상도 이미지 합성 연구" [1] 를 기반으로 하여, Stability AI와 Runway ML 등의 지원을 받아 개발된 딥러닝 인공지능 모델이다. 1. You can use DynamicPrompt Extantion with prompt like: {1-15$$__all__} to get completely random results. Utilizing the latent diffusion model, a variant of the diffusion model, it effectively removes even the strongest noise from data. Sensitive Content. a CompVis. Hな表情の呪文・プロンプト. Our codebase for the diffusion models builds heavily on OpenAI’s ADM codebase and Thanks for open-sourcing! CompVis initial stable diffusion release; Patrick’s implementation of the streamlit demo for inpainting. g. r/StableDiffusion. 管不了了. StableSwarmUI, A Modular Stable Diffusion Web-User-Interface, with an emphasis on making powertools easily accessible, high performance, and extensibility. It is primarily used to generate detailed images conditioned on text descriptions. Stable Diffusion Prompts. Free-form inpainting is the task of adding new content to an image in the regions specified by an arbitrary binary mask. 1. youtube. Controlnet v1. the command-line version of Stable Diffusion, you just add a full colon followed by a decimal number to the word you want to emphasize. Features. (Open in Colab) Build your own Stable Diffusion UNet model from scratch in a notebook. Next, make sure you have Pyhton 3. 34k. 8k stars Watchers. Fooocus is an image generating software (based on Gradio ). Stable Diffusion 🎨. It is trained on 512x512 images from a subset of the LAION-5B database. Load safetensors. People have asked about the models I use and I've promised to release them, so here they are. For Stable Diffusion, we started with the FP32 version 1-5 open-source model from Hugging Face and made optimizations through quantization, compilation, and hardware acceleration to run it on a phone powered by Snapdragon 8 Gen 2 Mobile Platform.